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Jonathan Granskog

Publications and source records attributed to Jonathan Granskog.

4 recordsLinked to original sources

Runway Characters: Serving a 57B Video Diffusion Transformer for Real-Time Conversation

Conversational video characters must generate video as fast as it plays, for one user at a time: at 24 fps, a sustained 166.7 ms per four-frame chunk at batch size one, with only a few chunks of buffering. We describe the system behind Runway Characters, which serves a $\approx$57B-parameter autoregressive video diffusion transformer (four denoising passes per chunk) at 704$\times$1088 and 24 fps on four H100 GPUs with a voice agent in the loop. Each chunk takes 240 attention calls, each reading a 5,788-slot KV cache in which only 757 tokens are new, and 29% of the original step was whole-cache work repeated in every call. We move that work to write time (a ring cache, RoPE and FP8 quantization applied once per write, and a delta rotation when positions reset), fuse the tensor-parallel all-reduce into the FP8 GEMM with in-switch reduction, and give the decoder that turns latents into video frames its own GPU. Without the fused kernel, our optimizations cut the diffusion step from 240.9 to 159.3 ms (1.51$\times$, 24.5 fps sustained) with no detected degradation in lip-sync scores; the fused kernel, rolled back once in production and since re-enabled, reaches 147.4 ms (1.63$\times$, 43% of FP8 peak). In seven weeks of production the worker averaged above 23.5 fps in 99.5% of sessions that streamed for at least a second, and our pipeline takes 0.49 s of a 2.1 s median conversational turn. On the same GPUs, the two open-source real-time video systems we benchmarked, at four denoising passes, miss real time or add over a second of latency (one by estimate), at a third to half of our resolution.

cs.DC↗

Dual-Process Image Generation

Prior methods for controlling image generation are limited in their ability to be taught new tasks. In contrast, vision-language models, or VLMs, can learn tasks in-context and produce the correct outputs for a given input. We propose a dual-process distillation scheme that allows feed-forward image generators to learn new tasks from deliberative VLMs. Our scheme uses a VLM to rate the generated images and backpropagates this gradient to update the weights of the image generator. Our general framework enables a wide variety of new control tasks through the same text-and-image based interface. We showcase a handful of applications of this technique for different types of control signals, such as commonsense inferences and visual prompts. With our method, users can implement multimodal controls for properties such as color palette, line weight, horizon position, and relative depth within a matter of minutes. Project page: https://dual-process.github.io.

cs.CV↗

Structure and Content-Guided Video Synthesis with Diffusion Models

Text-guided generative diffusion models unlock powerful image creation and editing tools. While these have been extended to video generation, current approaches that edit the content of existing footage while retaining structure require expensive re-training for every input or rely on error-prone propagation of image edits across frames. In this work, we present a structure and content-guided video diffusion model that edits videos based on visual or textual descriptions of the desired output. Conflicts between user-provided content edits and structure representations occur due to insufficient disentanglement between the two aspects. As a solution, we show that training on monocular depth estimates with varying levels of detail provides control over structure and content fidelity. Our model is trained jointly on images and videos which also exposes explicit control of temporal consistency through a novel guidance method. Our experiments demonstrate a wide variety of successes; fine-grained control over output characteristics, customization based on a few reference images, and a strong user preference towards results by our model.

cs.CV↗

CLIP-based Neural Neighbor Style Transfer for 3D Assets

We present a method for transferring the style from a set of images to a 3D object. The texture appearance of an asset is optimized with a differentiable renderer in a pipeline based on losses using pretrained deep neural networks. More specifically, we utilize a nearest-neighbor feature matching loss with CLIP-ResNet50 to extract the style from images. We show that a CLIP- based style loss provides a different appearance over a VGG-based loss by focusing more on texture over geometric shapes. Additionally, we extend the loss to support multiple images and enable loss-based control over the color palette combined with automatic color palette extraction from style images.

cs.GR↗